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Update app.py
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app.py
CHANGED
@@ -4,30 +4,35 @@ import numpy as np
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from diffusers import StableDiffusionImg2ImgPipeline
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from PIL import Image
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from typing import Generator, List
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device = "cuda" if torch.cuda.is_available() else "cpu"
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model_id = "nitrosocke/Ghibli-Diffusion"
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#
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pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
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model_id,
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torch_dtype=torch.
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)
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pipe = pipe.to(device)
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pipe.enable_attention_slicing()
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def resize_and_crop(image: Image.Image, target_size: int = 512) -> Image.Image:
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"""
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width, height = image.size
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image = image.crop((left, 0, right, height))
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elif height > width:
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top = (height - width) // 2
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bottom = top + width
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image = image.crop((0, top, width, bottom))
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return image.resize((target_size, target_size))
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def generate_ghibli_style(
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input_image: Image.Image,
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@@ -35,90 +40,103 @@ def generate_ghibli_style(
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strength: float = 0.6,
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guidance_scale: float = 7.5
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) -> Generator[Image.Image, None, None]:
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"""
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prompt = "ghibli style, detailed anime portrait, studio ghibli, anime artwork"
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negative_prompt = "blurry, low quality, sketch, cartoon, 3d, deformed, disfigured"
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# Preprocess
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input_image = resize_and_crop(input_image)
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init_image = input_image.convert("RGB")
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# Prepare latent variables
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init_latents = pipe.vae.encode(
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init_latents = pipe.vae.config.scaling_factor * init_latents
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#
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pipe.scheduler.set_timesteps(steps, device=device)
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timesteps = pipe.scheduler.timesteps[int(steps * strength):]
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noise = torch.randn_like(init_latents)
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latents = pipe.scheduler.add_noise(init_latents, noise, timesteps[:1])
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#
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text_inputs = pipe.tokenizer(
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prompt,
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padding="max_length",
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max_length=pipe.tokenizer.model_max_length,
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return_tensors="pt"
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)
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text_embeddings = pipe.text_encoder(text_inputs.input_ids.to(device))[0]
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# Unconditional embedding
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uncond_input = pipe.tokenizer(
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[negative_prompt]
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padding="max_length",
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max_length=text_embeddings.shape[1],
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return_tensors="pt"
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)
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uncond_embeddings = pipe.text_encoder(uncond_input.input_ids.to(device))[0]
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# Classifier-free guidance
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text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
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# Diffusion process
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for i, t in enumerate(gr.Progress().tqdm(timesteps, desc="Generating")):
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#
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# Decode and yield image
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with torch.no_grad():
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image = pipe.vae.decode(latents / pipe.vae.config.scaling_factor, return_dict=False)[0]
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image = pipe.image_processor.postprocess(image, output_type="pil")[0]
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yield image
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# Gradio interface
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with gr.Blocks() as demo:
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gr.Markdown("# ✨ Studio Ghibli Style Transformer ✨")
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gr.Markdown("Upload a portrait photo to transform it into a Studio Ghibli-style artwork!")
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with gr.Row():
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with gr.Column():
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input_image = gr.Image(label="Input Image", type="pil")
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steps_slider = gr.Slider(10,
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strength_slider = gr.Slider(0.
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generate_btn = gr.Button("✨ Transform!", variant="primary")
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with gr.Column():
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gallery = gr.Gallery(
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label="Generation Progress",
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show_label=True,
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columns=
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preview=True,
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object_fit="contain",
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height=600
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concurrency_limit=1
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)
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if __name__ == "__main__":
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demo.launch()
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from diffusers import StableDiffusionImg2ImgPipeline
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from PIL import Image
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from typing import Generator, List
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import gc
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import os
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# Configure CPU optimization
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os.environ["OMP_NUM_THREADS"] = "1"
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os.environ["MKL_NUM_THREADS"] = "1"
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torch.set_num_threads(1)
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device = "cuda" if torch.cuda.is_available() else "cpu"
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model_id = "nitrosocke/Ghibli-Diffusion"
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# Memory-optimized pipeline loading
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pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
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model_id,
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torch_dtype=torch.float32, # Keep float32 for CPU stability
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)
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pipe = pipe.to(device)
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pipe.enable_attention_slicing(slice_size=4)
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pipe.enable_sequential_cpu_offload() if device == "cuda" else None
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def resize_and_crop(image: Image.Image, target_size: int = 512) -> Image.Image:
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"""Optimized image preprocessing with downsampling"""
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width, height = image.size
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scale = max(target_size/width, target_size/height)
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image = image.resize((int(width*scale), int(height*scale)), Image.LANCZOS)
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width, height = image.size
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left = (width - target_size) // 2
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top = (height - target_size) // 2
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return image.crop((left, top, left+target_size, top+target_size))
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def generate_ghibli_style(
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input_image: Image.Image,
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strength: float = 0.6,
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guidance_scale: float = 7.5
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) -> Generator[Image.Image, None, None]:
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"""Memory-optimized generator with aggressive cleanup"""
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prompt = "ghibli style, detailed anime portrait, studio ghibli, anime artwork"
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negative_prompt = "blurry, low quality, sketch, cartoon, 3d, deformed, disfigured"
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# Preprocess with garbage collection
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input_image = resize_and_crop(input_image)
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init_image = input_image.convert("RGB")
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del input_image
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gc.collect()
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# Prepare latent variables with memory mapping
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init_tensor = pipe.image_processor.preprocess(init_image).to(device=device, dtype=torch.float32)
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init_latents = pipe.vae.encode(init_tensor).latent_dist.sample()
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init_latents = pipe.vae.config.scaling_factor * init_latents
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del init_tensor
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gc.collect()
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# Configure scheduler
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pipe.scheduler.set_timesteps(steps, device=device)
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timesteps = pipe.scheduler.timesteps[int(steps * strength):]
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noise = torch.randn_like(init_latents, device=device)
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latents = pipe.scheduler.add_noise(init_latents, noise, timesteps[:1])
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del init_latents, noise
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gc.collect()
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# Memory-efficient text encoding
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text_inputs = pipe.tokenizer(
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prompt,
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padding="max_length",
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max_length=pipe.tokenizer.model_max_length,
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return_tensors="pt"
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)
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text_embeddings = pipe.text_encoder(text_inputs.input_ids.to(device))[0].to(torch.float32)
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uncond_input = pipe.tokenizer(
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[negative_prompt],
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padding="max_length",
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max_length=text_embeddings.shape[1],
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return_tensors="pt"
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)
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uncond_embeddings = pipe.text_encoder(uncond_input.input_ids.to(device))[0].to(torch.float32)
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text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
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del uncond_embeddings, uncond_input, text_inputs
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gc.collect()
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# Diffusion process with memory cleanup
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for i, t in enumerate(gr.Progress().tqdm(timesteps, desc="Generating")):
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# Memory-optimized UNet inference
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with torch.inference_mode():
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latent_model_input = torch.cat([latents] * 2)
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latent_model_input = pipe.scheduler.scale_model_input(latent_model_input, t)
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noise_pred = pipe.unet(
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latent_model_input,
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t,
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encoder_hidden_states=text_embeddings,
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return_dict=False,
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)[0]
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noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
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noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
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latents = pipe.scheduler.step(noise_pred, t, latents).prev_sample
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# Memory-efficient decoding
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with torch.no_grad():
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image = pipe.vae.decode(latents / pipe.vae.config.scaling_factor, return_dict=False)[0]
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image = pipe.image_processor.postprocess(image, output_type="pil")[0]
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yield image
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# Aggressive memory cleanup
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del latent_model_input, noise_pred, noise_pred_uncond, noise_pred_text
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gc.collect()
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# Final cleanup
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del latents, text_embeddings
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gc.collect()
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# Gradio interface
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with gr.Blocks() as demo:
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gr.Markdown("# ✨ Studio Ghibli Style Transformer (CPU Optimized) ✨")
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gr.Markdown("Upload a portrait photo to transform it into a Studio Ghibli-style artwork (max 10GB RAM usage)!")
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with gr.Row():
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with gr.Column():
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input_image = gr.Image(label="Input Image", type="pil")
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steps_slider = gr.Slider(10, 40, value=25, step=5, label="Number of Steps")
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strength_slider = gr.Slider(0.4, 0.8, value=0.6, step=0.1, label="Transformation Strength")
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generate_btn = gr.Button("✨ Transform!", variant="primary")
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with gr.Column():
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gallery = gr.Gallery(
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label="Generation Progress",
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show_label=True,
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columns=4,
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preview=True,
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object_fit="contain",
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height=600
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concurrency_limit=1
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)
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if __name__ == "__main__":
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demo.queue(concurrency_count=1)
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demo.launch()
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